Custom Diffusion Model Text-to-Image Generator

A custom implementation of a text-to-image generation model usiing stable diffusion weights.

Features

• Custom diffusion model implementation
• Configurable generation parameters
• User-friendly Gradio interface
• Automatic model downloading

Usage

1. Enter your prompt in the text box.
2. (Optional) Adjust advanced settings:
   • Negative prompt
   • Seed
   • Image dimensions
   • Guidance scale
   • Number of inference steps
3. Click "Run" to generate the image.

Technical Details

• Model: Stable Diffusion v1.5 with custom modifications
• Architecture: Custom UNet
• Default settings:
  • Image size: 512x512
  • Guidance scale: 7.5
  • Inference steps: 50
  • Seed: 42 (randomizable)

Examples

Try these example prompts:

• "A ultra sharp photorealistic painting of a futuristic cityscape at night with neon lights and flying cars"
• "A serene mountain landscape at sunset with snow-capped peaks and a clear lake reflection"
• "A detailed portrait of a cyberpunk character with glowing neon implants and holographic tattoos"

Installation

The model will be automatically downloaded on first run. No manual installation required.

License

This project is available under the MIT License.

Try it out on Hugging Face Spaces

Check out the live demo on Hugging Face Spaces: torinriley/Diffusion